VAE HERE: <a target="_blank" rel="ugc" href="https://civarchive.com/models/146075">https://civarchive.com/models/146075</a>"SD-XL Inpainting 0.1 is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask.The SD-XL Inpainting 0.1 was initialized with the <code>stable-diffusion-xl-base-1.0</code> weights. The model is trained for 40k steps at resolution 1024x1024 and 5% dropping of the text-conditioning to improve classifier-free classifier-free guidance sampling. For inpainting, the UNet has 5 additional input channels (4 for the encoded masked-image and 1 for the mask itself) whose weights were zero-initialized after restoring the non-inpainting checkpoint. During training, we generate synthetic masks and, in 25% mask everything."Guide for CompfyUI: <a target="_blank" rel="ugc" href="https://mybyways.com/blog/using-the-sdxl-inpainting-01-model">https://mybyways.com/blog/using-the-sdxl-inpainting-01-model</a><a target="_blank" rel="ugc" href="https://huggingface.co/diffusers/stable-diffusion-xl-1.0-inpainting-0.1">source: https://huggingface.co/diffusers/stable-diffusion-xl-1.0-inpainting-0.1</a>

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